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The former one was generated using SDXL 0.9 vs 1.0. The new one has been generated with diffusers: import torch from diffusers import StableDiffusionXLPipeline, DDIMScheduler noise_scheduler = DDIMScheduler( num_train_timesteps=1000, beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", clip_sample=False, set_alpha_to_one=False, steps_offset=1, ) base_model_path = "/path/to/stabilityai/stable-diffusion-xl-base-1.0" device = "cuda" prompt = "a cute cat, detailed high-quality professional image" negative_prompt = "lowres, bad anatomy, bad hands, cropped, worst quality" seed = 2 pipe = StableDiffusionXLPipeline.from_pretrained(base_model_path, scheduler=noise_scheduler, torch_dtype=torch.float16, add_watermarker=False) pipe = pipe.to(device) generator = torch.Generator(device).manual_seed(seed) images = pipe(prompt=prompt, negative_prompt=negative_prompt, num_inference_steps=30, generator=generator).images |
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.. | ||
test_diffusion_ref | ||
__init__.py | ||
test_diffusion.py | ||
test_preprocessors.py |